![]() An image generated with Stable Diffusion based on the text prompt "a photograph of an astronaut riding a horse" | |
Original author(s) | Runway, CompVis, and Stability AI |
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Developer(s) | Stability AI |
Initial release | August 22, 2022 |
Stable release | SDXL 1.0 (model)[1]
/ July 26, 2023 |
Repository | |
Written in | Python[2] |
Type | Text-to-image model |
License | Creative ML OpenRAIL-M |
Website | stability![]() |
Stable Diffusion is a deep learning, text-to-image model released in 2022 based on diffusion techniques. The generative artificial intelligence technology is the premier product of Stability AI and is considered to be a part of the ongoing artificial intelligence boom.
It is primarily used to generate detailed images conditioned on text descriptions, though it can also be applied to other tasks such as inpainting, outpainting, and generating image-to-image translations guided by a text prompt.[3] Its development involved researchers from the CompVis Group at Ludwig Maximilian University of Munich and Runway with a computational donation from Stability and training data from non-profit organizations.[4][5][6][7]
Stable Diffusion is a latent diffusion model, a kind of deep generative artificial neural network. Its code and model weights have been released publicly,[8] and it can run on most consumer hardware equipped with a modest GPU with at least 4 GB VRAM. This marked a departure from previous proprietary text-to-image models such as DALL-E and Midjourney which were accessible only via cloud services.[9][10]
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